Stable Diffusion AI Image Generator
Stable Diffusion is an advanced AI image model developed by Stability AI, designed for generating high-quality images. Released in 2022, it utilizes a technique known as latent diffusion, which combines generative modeling with diffusion processes to create images that closely resemble real-world visuals. Try Stable Diffusion on Pollo AI!
Key Features of Stable Diffusion AI
- Text-to-Image Generation: Convert text prompts into coherent and visually appealing images
- Image-to-Image Generation: Use both a text prompt and an initial image to create new images
- Inpainting: Remove or replace objects in an image
- Outpainting:Extend existing images with new, contextually consistent content
Text-to-Image Generation
The model excels at converting text prompts into coherent and visually appealing images. Users can input descriptive phrases, and Stable Diffusion generates corresponding visuals that capture the essence of the text.
| Text prompt | Output image |
| scene of a giant ancient tortoise with a fantasy city built on its back. The tortoise's shell is covered in lush, dense forest with towering trees and a hidden, misty village nestled in the foliage. The city consists of intricately designed buildings that blend seamlessly with the natural environment, featuring rope bridges connecting different sections of the city. |
 |
| the four elements in a beautiful glass box within an intricate glass box within a gorgeous glass box within a glass box. Ethereal, elements! Hyper-detailed, intricate, masterpiece inside a glass box. |
 |
| photo of three potions: the first potion is blue with the label "MANA", the second potion is red with the label "HEALTH", the third potion is green with the label "POISON". Old apothecary. |
 |
| collage art 'We're Leaving For the Future' 1980s #vaporwave aesthetic internet art glued layered magazine cutout image shape scrap, torn ragged paper art, BASIC code, halftone, #pixelart. |
 |
Image to Image Generation
Stable Diffusion AI's image-to-image generation uses both a text prompt and an initial image to create new images that share characteristics with the original. Instead of starting from random noise, the model adds noise to the initial image and then denoises it based on the text prompt. This allows the model to capture general features from the initial image.
| Input image | Prompt | Output image |
 |
A rainbow coloured tiger |
 |
Inpainting
Stable Diffusion AI's Inpainting feature is a powerful tool for editing images. It allows you to remove objects from an image or replace an object with another. It ensures the content filled in is seamless and natural.
| Inputs | Prompt | Output image |
Image:
 Mask:  |
An orange cat sitting on a bench |
 |
Outpainting
Stable Diffusion's Outpainting feature allows users to extend existing images beyond their original borders, creating new, contextually consistent content. It uses AI to generate new pixels that seamlessly expand the image's boundaries.
| Inputs | Output image |
 |
 |
 |
 |
How to Use Stable Diffusion on Pollo AI
Select the Stable Diffusion Model
Go to the Pollo AI AI image generator and choose the Stable Diffusion model.
Enter Your Prompt
Enter your text prompt and select other customizable options.
Create Image
Click Create to start the image generation process.
YouTube Videos About Stable Diffusion
What People Are Talking About Stable Diffusion on Reddit
What People Are Talking About Stable Diffusion on X
Some thoughts on Stable Diffusion 3 medium #SD3
— Emad (@EMostaque) June 14, 2024
1. It’s a good model with a blend of speed & performance
2. It was iteratively trained by Robin’s team & rest of Stability AI team to blend wide use but also be good out of the box
3. It’s clear some of the safety alignment…
Just tested Stable Diffusion 3.5 Large locally in Comfy UI and trust me, if you've been using FLUX, this is a MAJOR step back. It's got average image quality + bad hands/anatomy. It only excels at doing different styles. Just stick with Flux. Honest opinion. #AI pic.twitter.com/XkYiw3h8wi
— Travis Davids (@MrDavids1) October 23, 2024
In July 2022, Stable Diffusion hadn't been released and I was playing with an early access version of it.
— TomLikesRobots🤖 (@TomLikesRobots) January 22, 2025
It blows me away that just two and a half years later I can use the same prompt to generate a near life like video clip.
"A white horse on a Black Sand Beach in Iceland". https://t.co/Tn3Yricbwq pic.twitter.com/lkqB9KfsVr
I feel very conflicted about the Stable Diffusion open source release.
— Joshua Achiam (@jachiam0) September 10, 2022
Stable Diffusion 3 model is a novel approach that combines the concepts of diffusion with flow matching and timestamp sampling.
— Satya Mallick (@LearnOpenCV) November 19, 2024
This week's blog post at LearnOpenCV will provide you an intuitive approach on Stable Diffusion 3 and 3.5 which can serve as a stepping stone for those… pic.twitter.com/1H3J8MU9cg
🚀 Customers can now access Stable Diffusion 3.5 Large in Amazon SageMaker JumpStart: https://t.co/XGDAyapWZW
— Swami Sivasubramanian (@SwamiSivasubram) November 14, 2024
At 8.1B parameters, Stable Diffusion 3.5 Large is the most powerful text-to-image model in the Stable Diffusion family with superior quality & prompt adherence. For…
Stable Diffusion (via Midjourney) was my "aha" moment in AI three (?) years ago.
— emozilla (@theemozilla) February 12, 2025
I was AI adjacent working in automotive, but that was what immediately made me (almost literally) drop everything and focus on learning it (see https://t.co/R2KxmFyRHi)
Today my wife asked me to… pic.twitter.com/MmeifeDuE9
Let's take a look at the Stable Diffusion 3.5 variants: Stable Diffusion 3.5 Large, Stable Diffusion 3.5 Large Turbo, and Stable Diffusion 3.5 Medium. https://t.co/mPwrEQHbNT
— vast.ai (@vast_ai) December 16, 2024
Stable Diffusion 3 is announced
— Holly Herndon (@hollyherndon) February 23, 2024
Stability used @spawning_ Do Not Train registry, which has over 1.5B opt-out requests, to filter their datasets before training
Many more models will be released this year honoring opt-outs. I hope we are getting closer to it being standard! https://t.co/NRrnHJPgE0
Taking a look at people testing out Stable Diffusion 3 and tbh this goes hard. pic.twitter.com/Ii7bHqmX9Y
— Max Woolf (@minimaxir) June 12, 2024
Stable Diffusion on @daytonaio infra 🔥 https://t.co/EgD8htZDJp
— Ivan Burazin (@ivanburazin) December 13, 2024
Discover Other AI Image Models
FAQs
What does Stable Diffusion do?
Stable Diffusion is a deep learning, text-to-image model that generates detailed, realistic images from text descriptions. It is a generative AI model based on diffusion techniques, primarily used to create original images from text prompts. Stable Diffusion can also perform tasks such as inpainting, outpainting, and image-to-image translations guided by a text prompt.
How do I write good prompts for Stable Diffusion?
Include details about the subject, style, composition, lighting, and any other relevant attributes. And also specify what you don't want in the image to prevent unwanted artifacts or styles.
Can I use Stable Diffusion for free?
Yes, the model itself is offered for free and you can run it through cloud-based services or locally on your machine.
Is Stable Diffusion easy to use for beginners?
Stable Diffusion may have a learning curve for beginners, but you can try it on Pollo AI. We provide an intuitive interface that makes Stable Diffusion image generation accessible to both professionals and amateurs.
